[docs] Clean up pipeline apis (#3905)

* start with stable diffusion * fix * finish stable diffusion pipelines * fix path to pipeline output * fix flax paths * fix copies * add up to score sde ve * finish first pass of pipelines * fix copies * second review * align doc titles * more review fixes * final review

[docs] Clean up pipeline apis (#3905)
* start with stable diffusion * fix * finish stable diffusion pipelines * fix path to pipeline output * fix flax paths * fix copies * add up to score sde ve * finish first pass of pipelines * fix copies * second review * align doc titles * more review fixes * final review
a69754bb · Steven Liu · GitHub · bcc570b9 · a69754bb · a69754bb
Unverified Commit a69754bb authored Jul 21, 2023 by Steven Liu Committed by GitHub Jul 21, 2023
20 changed files
--- a/docs/source/en/_toctree.yml
+++ b/docs/source/en/_toctree.yml
@@ -181,7 +181,7 @@
    - local: api/pipelines/alt_diffusion
      title: AltDiffusion
    - local: api/pipelines/attend_and_excite
-      title: Attend and Excite
+      title: Attend-and-Excite
    - local: api/pipelines/audio_diffusion
      title: Audio Diffusion
    - local: api/pipelines/audioldm
@@ -211,7 +211,7 @@
    - local: api/pipelines/latent_diffusion
      title: Latent Diffusion
    - local: api/pipelines/panorama
-      title: MultiDiffusion Panorama
+      title: MultiDiffusion
    - local: api/pipelines/paint_by_example
      title: PaintByExample
    - local: api/pipelines/paradigms
@@ -236,15 +236,15 @@
      - local: api/pipelines/stable_diffusion/overview
        title: Overview
      - local: api/pipelines/stable_diffusion/text2img
-        title: Text-to-Image
+        title: Text-to-image
      - local: api/pipelines/stable_diffusion/img2img
-        title: Image-to-Image
+        title: Image-to-image
      - local: api/pipelines/stable_diffusion/inpaint
-        title: Inpaint
+        title: Inpainting
      - local: api/pipelines/stable_diffusion/depth2img
-        title: Depth-to-Image
+        title: Depth-to-image
      - local: api/pipelines/stable_diffusion/image_variation
-        title: Image-Variation
+        title: Image variation
      - local: api/pipelines/stable_diffusion/stable_diffusion_safe
        title: Safe Stable Diffusion
      - local: api/pipelines/stable_diffusion/stable_diffusion_2
@@ -252,9 +252,9 @@
      - local: api/pipelines/stable_diffusion/stable_diffusion_xl
        title: Stable Diffusion XL
      - local: api/pipelines/stable_diffusion/latent_upscale
-        title: Stable-Diffusion-Latent-Upscaler
+        title: Latent upscaler
      - local: api/pipelines/stable_diffusion/upscale
-        title: Super-Resolution
+        title: Super-resolution
      - local: api/pipelines/stable_diffusion/ldm3d_diffusion
        title: LDM3D Text-to-(RGB, Depth)
      - local: api/pipelines/stable_diffusion/adapter
@@ -265,11 +265,11 @@
    - local: api/pipelines/stochastic_karras_ve
      title: Stochastic Karras VE
    - local: api/pipelines/model_editing
-      title: Text-to-Image Model Editing
+      title: Text-to-image model editing
    - local: api/pipelines/text_to_video
-      title: Text-to-Video
+      title: Text-to-video
    - local: api/pipelines/text_to_video_zero
-      title: Text-to-Video Zero
+      title: Text2Video-Zero
    - local: api/pipelines/unclip
      title: UnCLIP
    - local: api/pipelines/latent_diffusion_uncond

--- a/docs/source/en/api/pipelines/alt_diffusion.mdx
+++ b/docs/source/en/api/pipelines/alt_diffusion.mdx
@@ -12,72 +12,36 @@ specific language governing permissions and limitations under the License.
 # AltDiffusion
-AltDiffusion was proposed in [AltCLIP: Altering the Language Encoder in CLIP for Extended Language Capabilities](https://arxiv.org/abs/2211.06679) by Zhongzhi Chen, Guang Liu, Bo-Wen Zhang, Fulong Ye, Qinghong Yang, Ledell Wu.
+AltDiffusion was proposed in [AltCLIP: Altering the Language Encoder in CLIP for Extended Language Capabilities](https://huggingface.co/papers/2211.06679) by Zhongzhi Chen, Guang Liu, Bo-Wen Zhang, Fulong Ye, Qinghong Yang, Ledell Wu.
-The abstract of the paper is the following:
+The abstract from the paper is:
 *In this work, we present a conceptually simple and effective method to train a strong bilingual multimodal representation model. Starting from the pretrained multimodal representation model CLIP released by OpenAI, we switched its text encoder with a pretrained multilingual text encoder XLM-R, and aligned both languages and image representations by a two-stage training schema consisting of teacher learning and contrastive learning. We validate our method through evaluations of a wide range of tasks. We set new state-of-the-art performances on a bunch of tasks including ImageNet-CN, Flicker30k- CN, and COCO-CN. Further, we obtain very close performances with CLIP on almost all tasks, suggesting that one can simply alter the text encoder in CLIP for extended capabilities such as multilingual understanding.*
-*Overview*:
-| Pipeline | Tasks | Colab | Demo
-|---|---|:---:|:---:|
-| [pipeline_alt_diffusion.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/alt_diffusion/pipeline_alt_diffusion.py) | *Text-to-Image Generation* | - | -
-| [pipeline_alt_diffusion_img2img.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/alt_diffusion/pipeline_alt_diffusion_img2img.py) | *Image-to-Image Text-Guided Generation* | - |-
 ## Tips
- AltDiffusion is conceptually exactly the same as [Stable Diffusion](./stable_diffusion/overview).
+`AltDiffusion` is conceptually the same as [Stable Diffusion](./stable_diffusion/overview).
- *Run AltDiffusion*
-AltDiffusion can be tested very easily with the [`AltDiffusionPipeline`], [`AltDiffusionImg2ImgPipeline`] and the `"BAAI/AltDiffusion-m9"` checkpoint exactly in the same way it is shown in the [Conditional Image Generation Guide](../../using-diffusers/conditional_image_generation) and the [Image-to-Image Generation Guide](../../using-diffusers/img2img).
- *How to load and use different schedulers.*
-The alt diffusion pipeline uses [`DDIMScheduler`] scheduler by default. But `diffusers` provides many other schedulers that can be used with the alt diffusion pipeline such as [`PNDMScheduler`], [`LMSDiscreteScheduler`], [`EulerDiscreteScheduler`], [`EulerAncestralDiscreteScheduler`] etc.
-To use a different scheduler, you can either change it via the [`ConfigMixin.from_config`] method or pass the `scheduler` argument to the `from_pretrained` method of the pipeline. For example, to use the [`EulerDiscreteScheduler`], you can do the following:
-```python
->>> from diffusers import AltDiffusionPipeline, EulerDiscreteScheduler
->>> pipeline = AltDiffusionPipeline.from_pretrained("BAAI/AltDiffusion-m9")
->>> pipeline.scheduler = EulerDiscreteScheduler.from_config(pipeline.scheduler.config)
->>> # or
->>> euler_scheduler = EulerDiscreteScheduler.from_pretrained("BAAI/AltDiffusion-m9", subfolder="scheduler")
->>> pipeline = AltDiffusionPipeline.from_pretrained("BAAI/AltDiffusion-m9", scheduler=euler_scheduler)
-```
- *How to convert all use cases with multiple or single pipeline*
-If you want to use all possible use cases in a single `DiffusionPipeline` we recommend using the `components` functionality to instantiate all components in the most memory-efficient way:
-```python
+<Tip>
->>> from diffusers import (
-...     AltDiffusionPipeline,
-...     AltDiffusionImg2ImgPipeline,
-... )
->>> text2img = AltDiffusionPipeline.from_pretrained("BAAI/AltDiffusion-m9")
+Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
->>> img2img = AltDiffusionImg2ImgPipeline(**text2img.components)
->>> # now you can use text2img(...) and img2img(...) just like the call methods of each respective pipeline
+</Tip>
-```
-## AltDiffusionPipelineOutput
-[[autodoc]] pipelines.alt_diffusion.AltDiffusionPipelineOutput
-	- all
-	- __call__
 ## AltDiffusionPipeline
 [[autodoc]] AltDiffusionPipeline
 	- all
 	- __call__
 ## AltDiffusionImg2ImgPipeline
 [[autodoc]] AltDiffusionImg2ImgPipeline
 	- all
 	- __call__
+## AltDiffusionPipelineOutput
+[[autodoc]] pipelines.alt_diffusion.AltDiffusionPipelineOutput
+	- all
+	- __call__
\ No newline at end of file
--- a/docs/source/en/api/pipelines/attend_and_excite.mdx
+++ b/docs/source/en/api/pipelines/attend_and_excite.mdx
@@ -10,66 +10,28 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
 specific language governing permissions and limitations under the License.
 -->
-# Attend and Excite: Attention-Based Semantic Guidance for Text-to-Image Diffusion Models
+# Attend-and-Excite
-## Overview
+Attend-and-Excite for Stable Diffusion was proposed in [Attend-and-Excite: Attention-Based Semantic Guidance for Text-to-Image Diffusion Models](https://attendandexcite.github.io/Attend-and-Excite/) and provides textual attention control over image generation.
-Attend and Excite for Stable Diffusion was proposed in [Attend-and-Excite: Attention-Based Semantic Guidance for Text-to-Image Diffusion Models](https://attendandexcite.github.io/Attend-and-Excite/) and provides textual attention control over the image generation.
+The abstract from the paper is:
-The abstract of the paper is the following:
 *Text-to-image diffusion models have recently received a lot of interest for their astonishing ability to produce high-fidelity images from text only. However, achieving one-shot generation that aligns with the user's intent is nearly impossible, yet small changes to the input prompt often result in very different images. This leaves the user with little semantic control. To put the user in control, we show how to interact with the diffusion process to flexibly steer it along semantic directions. This semantic guidance (SEGA) allows for subtle and extensive edits, changes in composition and style, as well as optimizing the overall artistic conception. We demonstrate SEGA's effectiveness on a variety of tasks and provide evidence for its versatility and flexibility.*
-Resources
+You can find additional information about Attend-and-Excite on the [project page](https://attendandexcite.github.io/Attend-and-Excite/), the [original codebase](https://github.com/AttendAndExcite/Attend-and-Excite), or try it out in a [demo](https://huggingface.co/spaces/AttendAndExcite/Attend-and-Excite).
-* [Project Page](https://attendandexcite.github.io/Attend-and-Excite/)
-* [Paper](https://arxiv.org/abs/2301.13826)
-* [Original Code](https://github.com/AttendAndExcite/Attend-and-Excite)
-* [Demo](https://huggingface.co/spaces/AttendAndExcite/Attend-and-Excite)
-## Available Pipelines:
-| Pipeline | Tasks | Colab | Demo
-|---|---|:---:|:---:|
-| [pipeline_semantic_stable_diffusion_attend_and_excite.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_semantic_stable_diffusion_attend_and_excite) | *Text-to-Image Generation* | - | https://huggingface.co/spaces/AttendAndExcite/Attend-and-Excite
-### Usage example
-```python
+<Tip>
-import torch
-from diffusers import StableDiffusionAttendAndExcitePipeline
-model_id = "CompVis/stable-diffusion-v1-4"
+Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
-pipe = StableDiffusionAttendAndExcitePipeline.from_pretrained(model_id, torch_dtype=torch.float16).to("cuda")
-pipe = pipe.to("cuda")
-prompt = "a cat and a frog"
-# use get_indices function to find out indices of the tokens you want to alter
-pipe.get_indices(prompt)
-token_indices = [2, 5]
-seed = 6141
-generator = torch.Generator("cuda").manual_seed(seed)
-images = pipe(
-    prompt=prompt,
-    token_indices=token_indices,
-    guidance_scale=7.5,
-    generator=generator,
-    num_inference_steps=50,
-    max_iter_to_alter=25,
-).images
-image = images[0]
-image.save(f"../images/{prompt}_{seed}.png")
-```
+</Tip>
 ## StableDiffusionAttendAndExcitePipeline
 [[autodoc]] StableDiffusionAttendAndExcitePipeline
 	- all
 	- __call__
+## StableDiffusionPipelineOutput
+[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput
\ No newline at end of file
--- a/docs/source/en/api/pipelines/audio_diffusion.mdx
+++ b/docs/source/en/api/pipelines/audio_diffusion.mdx
@@ -12,87 +12,26 @@ specific language governing permissions and limitations under the License.
 # Audio Diffusion
-## Overview
+[Audio Diffusion](https://github.com/teticio/audio-diffusion) is by Robert Dargavel Smith, and it leverages the recent advances in image generation from diffusion models by converting audio samples to and from Mel spectrogram images.
-[Audio Diffusion](https://github.com/teticio/audio-diffusion) by Robert Dargavel Smith.
+The original codebase, training scripts and example notebooks can be found at [teticio/audio-diffusion](https://github.com/teticio/audio-diffusion).
-Audio Diffusion leverages the recent advances in image generation using diffusion models by converting audio samples to
+<Tip>
-and from mel spectrogram images.
-The original codebase of this implementation can be found [here](https://github.com/teticio/audio-diffusion), including
+Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
-training scripts and example notebooks.
-## Available Pipelines:
+</Tip>
-| Pipeline | Tasks | Colab
-|---|---|:---:|
-| [pipeline_audio_diffusion.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/audio_diffusion/pipeline_audio_diffusion.py) | *Unconditional Audio Generation* | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/teticio/audio-diffusion/blob/master/notebooks/audio_diffusion_pipeline.ipynb) |
-## Examples:
-### Audio Diffusion
-```python
-import torch
-from IPython.display import Audio
-from diffusers import DiffusionPipeline
-device = "cuda" if torch.cuda.is_available() else "cpu"
-pipe = DiffusionPipeline.from_pretrained("teticio/audio-diffusion-256").to(device)
-output = pipe()
-display(output.images[0])
-display(Audio(output.audios[0], rate=pipe.mel.get_sample_rate()))
-```
-### Latent Audio Diffusion
-```python
-import torch
-from IPython.display import Audio
-from diffusers import DiffusionPipeline
-device = "cuda" if torch.cuda.is_available() else "cpu"
-pipe = DiffusionPipeline.from_pretrained("teticio/latent-audio-diffusion-256").to(device)
-output = pipe()
-display(output.images[0])
-display(Audio(output.audios[0], rate=pipe.mel.get_sample_rate()))
-```
-### Audio Diffusion with DDIM (faster)
-```python
-import torch
-from IPython.display import Audio
-from diffusers import DiffusionPipeline
-device = "cuda" if torch.cuda.is_available() else "cpu"
-pipe = DiffusionPipeline.from_pretrained("teticio/audio-diffusion-ddim-256").to(device)
-output = pipe()
-display(output.images[0])
-display(Audio(output.audios[0], rate=pipe.mel.get_sample_rate()))
-```
-### Variations, in-painting, out-painting etc.
-```python
-output = pipe(
-    raw_audio=output.audios[0, 0],
-    start_step=int(pipe.get_default_steps() / 2),
-    mask_start_secs=1,
-    mask_end_secs=1,
-)
-display(output.images[0])
-display(Audio(output.audios[0], rate=pipe.mel.get_sample_rate()))
-```
 ## AudioDiffusionPipeline
 [[autodoc]] AudioDiffusionPipeline
 	- all
 	- __call__
+## AudioPipelineOutput
+[[autodoc]] pipelines.AudioPipelineOutput
+## ImagePipelineOutput
+[[autodoc]] pipelines.ImagePipelineOutput
 ## Mel
 [[autodoc]] Mel
--- a/docs/source/en/api/pipelines/audioldm.mdx
+++ b/docs/source/en/api/pipelines/audioldm.mdx
@@ -12,73 +12,40 @@ specific language governing permissions and limitations under the License.
 # AudioLDM
-## Overview
+AudioLDM was proposed in [AudioLDM: Text-to-Audio Generation with Latent Diffusion Models](https://huggingface.co/papers/2301.12503) by Haohe Liu et al. Inspired by [Stable Diffusion](https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion/overview), AudioLDM
-AudioLDM was proposed in [AudioLDM: Text-to-Audio Generation with Latent Diffusion Models](https://arxiv.org/abs/2301.12503) by Haohe Liu et al.
-Inspired by [Stable Diffusion](https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion/overview), AudioLDM
 is a text-to-audio _latent diffusion model (LDM)_ that learns continuous audio representations from [CLAP](https://huggingface.co/docs/transformers/main/model_doc/clap)
 latents. AudioLDM takes a text prompt as input and predicts the corresponding audio. It can generate text-conditional
 sound effects, human speech and music.
-This pipeline was contributed by [sanchit-gandhi](https://huggingface.co/sanchit-gandhi). The original codebase can be found [here](https://github.com/haoheliu/AudioLDM).
+The abstract from the paper is:
-## Text-to-Audio
-The [`AudioLDMPipeline`] can be used to load pre-trained weights from [cvssp/audioldm-s-full-v2](https://huggingface.co/cvssp/audioldm-s-full-v2) and generate text-conditional audio outputs:
+*Text-to-audio (TTA) system has recently gained attention for its ability to synthesize general audio based on text descriptions. However, previous studies in TTA have limited generation quality with high computational costs. In this study, we propose AudioLDM, a TTA system that is built on a latent space to learn the continuous audio representations from contrastive language-audio pretraining (CLAP) latents. The pretrained CLAP models enable us to train LDMs with audio embedding while providing text embedding as a condition during sampling. By learning the latent representations of audio signals and their compositions without modeling the cross-modal relationship, AudioLDM is advantageous in both generation quality and computational efficiency. Trained on AudioCaps with a single GPU, AudioLDM achieves state-of-the-art TTA performance measured by both objective and subjective metrics (e.g., frechet distance). Moreover, AudioLDM is the first TTA system that enables various text-guided audio manipulations (e.g., style transfer) in a zero-shot fashion. Our implementation and demos are available at https://audioldm.github.io.*
-```python
+The original codebase can be found at [haoheliu/AudioLDM](https://github.com/haoheliu/AudioLDM). 
-from diffusers import AudioLDMPipeline
-import torch
-import scipy
-repo_id = "cvssp/audioldm-s-full-v2"
+## Tips
-pipe = AudioLDMPipeline.from_pretrained(repo_id, torch_dtype=torch.float16)
-pipe = pipe.to("cuda")
-prompt = "Techno music with a strong, upbeat tempo and high melodic riffs"
+When constructing a prompt, keep in mind:
-audio = pipe(prompt, num_inference_steps=10, audio_length_in_s=5.0).audios[0]
-# save the audio sample as a .wav file
+* Descriptive prompt inputs work best; you can use adjectives to describe the sound (for example, "high quality" or "clear") and make the prompt context specific (for example, "water stream in a forest" instead of "stream").
-scipy.io.wavfile.write("techno.wav", rate=16000, data=audio)
+* It's best to use general terms like "cat" or "dog" instead of specific names or abstract objects the model may not be familiar with.
-```
-### Tips
+During inference:
-Prompts:
+* The _quality_ of the predicted audio sample can be controlled by the `num_inference_steps` argument; higher steps give higher quality audio at the expense of slower inference.
-* Descriptive prompt inputs work best: you can use adjectives to describe the sound (e.g. "high quality" or "clear") and make the prompt context specific (e.g., "water stream in a forest" instead of "stream").
-* It's best to use general terms like 'cat' or 'dog' instead of specific names or abstract objects that the model may not be familiar with.
-Inference:
-* The _quality_ of the predicted audio sample can be controlled by the `num_inference_steps` argument: higher steps give higher quality audio at the expense of slower inference.
 * The _length_ of the predicted audio sample can be controlled by varying the `audio_length_in_s` argument.
-### How to load and use different schedulers
+<Tip>
-The AudioLDM pipeline uses [`DDIMScheduler`] scheduler by default. But `diffusers` provides many other schedulers
-that can be used with the AudioLDM pipeline such as [`PNDMScheduler`], [`LMSDiscreteScheduler`], [`EulerDiscreteScheduler`],
-[`EulerAncestralDiscreteScheduler`] etc. We recommend using the [`DPMSolverMultistepScheduler`] as it's currently the fastest
-scheduler there is.
-To use a different scheduler, you can either change it via the [`ConfigMixin.from_config`]
+Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
-method, or pass the `scheduler` argument to the `from_pretrained` method of the pipeline. For example, to use the
-[`DPMSolverMultistepScheduler`], you can do the following:
-```python
+</Tip>
->>> from diffusers import AudioLDMPipeline, DPMSolverMultistepScheduler
->>> import torch
->>> pipeline = AudioLDMPipeline.from_pretrained("cvssp/audioldm-s-full-v2", torch_dtype=torch.float16)
->>> pipeline.scheduler = DPMSolverMultistepScheduler.from_config(pipeline.scheduler.config)
->>> # or
->>> dpm_scheduler = DPMSolverMultistepScheduler.from_pretrained("cvssp/audioldm-s-full-v2", subfolder="scheduler")
->>> pipeline = AudioLDMPipeline.from_pretrained(
-...     "cvssp/audioldm-s-full-v2", scheduler=dpm_scheduler, torch_dtype=torch.float16
-... )
-```
 ## AudioLDMPipeline
 [[autodoc]] AudioLDMPipeline
 	- all
 	- __call__
+## StableDiffusionPipelineOutput
+[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput
\ No newline at end of file
--- a/docs/source/en/api/pipelines/consistency_models.mdx
+++ b/docs/source/en/api/pipelines/consistency_models.mdx
 # Consistency Models
-Consistency Models were proposed in [Consistency Models](https://arxiv.org/abs/2303.01469) by Yang Song, Prafulla Dhariwal, Mark Chen, and Ilya Sutskever.
+Consistency Models were proposed in [Consistency Models](https://huggingface.co/papers/2303.01469) by Yang Song, Prafulla Dhariwal, Mark Chen, and Ilya Sutskever.
-The abstract of the [paper](https://arxiv.org/pdf/2303.01469.pdf) is as follows:
+The abstract from the paper is:
 *Diffusion models have significantly advanced the fields of image, audio, and video generation, but they depend on an iterative sampling process that causes slow generation. To overcome this limitation, we propose consistency models, a new family of models that generate high quality samples by directly mapping noise to data. They support fast one-step generation by design, while still allowing multistep sampling to trade compute for sample quality. They also support zero-shot data editing, such as image inpainting, colorization, and super-resolution, without requiring explicit training on these tasks. Consistency models can be trained either by distilling pre-trained diffusion models, or as standalone generative models altogether. Through extensive experiments, we demonstrate that they outperform existing distillation techniques for diffusion models in one- and few-step sampling, achieving the new state-of-the-art FID of 3.55 on CIFAR-10 and 6.20 on ImageNet 64x64 for one-step generation. When trained in isolation, consistency models become a new family of generative models that can outperform existing one-step, non-adversarial generative models on standard benchmarks such as CIFAR-10, ImageNet 64x64 and LSUN 256x256. *
-Resources:
+The original codebase can be found at [openai/consistency_models](https://github.com/openai/consistency_models), and additional checkpoints are available at [openai](https://huggingface.co/openai).
-* [Paper](https://arxiv.org/abs/2303.01469)
+The pipeline was contributed by [dg845](https://github.com/dg845) and [ayushtues](https://huggingface.co/ayushtues). ❤️
-* [Original Code](https://github.com/openai/consistency_models)
-Available Checkpoints are:
+## Tips
- *cd_imagenet64_l2 (64x64 resolution)* [openai/consistency-model-pipelines](https://huggingface.co/openai/diffusers-cd_imagenet64_l2)
- *cd_imagenet64_lpips (64x64 resolution)* [openai/diffusers-cd_imagenet64_lpips](https://huggingface.co/openai/diffusers-cd_imagenet64_lpips)
- *ct_imagenet64 (64x64 resolution)* [openai/diffusers-ct_imagenet64](https://huggingface.co/openai/diffusers-ct_imagenet64)
- *cd_bedroom256_l2 (256x256 resolution)* [openai/diffusers-cd_bedroom256_l2](https://huggingface.co/openai/diffusers-cd_bedroom256_l2)
- *cd_bedroom256_lpips (256x256 resolution)* [openai/diffusers-cd_bedroom256_lpips](https://huggingface.co/openai/diffusers-cd_bedroom256_lpips)
- *ct_bedroom256 (256x256 resolution)* [openai/diffusers-ct_bedroom256](https://huggingface.co/openai/diffusers-ct_bedroom256)
- *cd_cat256_l2 (256x256 resolution)* [openai/diffusers-cd_cat256_l2](https://huggingface.co/openai/diffusers-cd_cat256_l2)
- *cd_cat256_lpips (256x256 resolution)* [openai/diffusers-cd_cat256_lpips](https://huggingface.co/openai/diffusers-cd_cat256_lpips)
- *ct_cat256 (256x256 resolution)* [openai/diffusers-ct_cat256](https://huggingface.co/openai/diffusers-ct_cat256)
-## Available Pipelines
+For an additional speed-up, use `torch.compile` to generate multiple images in <1 second:
-| Pipeline | Tasks | Demo | Colab |
+```diff
-|:---:|:---:|:---:|:---:|
+  import torch
-| [ConsistencyModelPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/pipeline_consistency_models.py) | *Unconditional Image Generation* | | |
+  from diffusers import ConsistencyModelPipeline
-This pipeline was contributed by our community members [dg845](https://github.com/dg845) and [ayushtues](https://huggingface.co/ayushtues) ❤️
+  device = "cuda"
+  # Load the cd_bedroom256_lpips checkpoint.
+  model_id_or_path = "openai/diffusers-cd_bedroom256_lpips"
+  pipe = ConsistencyModelPipeline.from_pretrained(model_id_or_path, torch_dtype=torch.float16)
+  pipe.to(device)
-## Usage Example
+ pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)
-```python
+  # Multistep sampling
-import torch
+  # Timesteps can be explicitly specified; the particular timesteps below are from the original Github repo:
+  # https://github.com/openai/consistency_models/blob/main/scripts/launch.sh#L83
-from diffusers import ConsistencyModelPipeline
+  for _ in range(10):
+      image = pipe(timesteps=[17, 0]).images[0]
-device = "cuda"
+      image.show()
-# Load the cd_imagenet64_l2 checkpoint.
-model_id_or_path = "openai/diffusers-cd_imagenet64_l2"
-pipe = ConsistencyModelPipeline.from_pretrained(model_id_or_path, torch_dtype=torch.float16)
-pipe.to(device)
-# Onestep Sampling
-image = pipe(num_inference_steps=1).images[0]
-image.save("consistency_model_onestep_sample.png")
-# Onestep sampling, class-conditional image generation
-# ImageNet-64 class label 145 corresponds to king penguins
-image = pipe(num_inference_steps=1, class_labels=145).images[0]
-image.save("consistency_model_onestep_sample_penguin.png")
-# Multistep sampling, class-conditional image generation
-# Timesteps can be explicitly specified; the particular timesteps below are from the original Github repo.
-# https://github.com/openai/consistency_models/blob/main/scripts/launch.sh#L77
-image = pipe(timesteps=[22, 0], class_labels=145).images[0]
-image.save("consistency_model_multistep_sample_penguin.png")
-```
-For an additional speed-up, one can also make use of `torch.compile`. Multiple images can be generated in <1 second as follows:
-```py
-import torch
-from diffusers import ConsistencyModelPipeline
-device = "cuda"
-# Load the cd_bedroom256_lpips checkpoint.
-model_id_or_path = "openai/diffusers-cd_bedroom256_lpips"
-pipe = ConsistencyModelPipeline.from_pretrained(model_id_or_path, torch_dtype=torch.float16)
-pipe.to(device)
-pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)
-# Multistep sampling
-# Timesteps can be explicitly specified; the particular timesteps below are from the original Github repo:
-# https://github.com/openai/consistency_models/blob/main/scripts/launch.sh#L83
-for _ in range(10):
-    image = pipe(timesteps=[17, 0]).images[0]
-    image.show()
 ```
 ## ConsistencyModelPipeline
 [[autodoc]] ConsistencyModelPipeline
    - all
    - __call__
+## ImagePipelineOutput
+[[autodoc]] pipelines.ImagePipelineOutput
\ No newline at end of file
--- a/docs/source/en/api/pipelines/controlnet.mdx
+++ b/docs/source/en/api/pipelines/controlnet.mdx
@@ -10,32 +10,19 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
 specific language governing permissions and limitations under the License.
 -->
-# Text-to-Image Generation with ControlNet Conditioning
+# ControlNet
-## Overview
+[Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang and Maneesh Agrawala.
-[Adding Conditional Control to Text-to-Image Diffusion Models](https://arxiv.org/abs/2302.05543) by Lvmin Zhang and Maneesh Agrawala.
+Using a pretrained model, we can provide control images (for example, a depth map) to control Stable Diffusion text-to-image generation so that it follows the structure of the depth image and fills in the details.
-Using the pretrained models we can provide control images (for example, a depth map) to control Stable Diffusion text-to-image generation so that it follows the structure of the depth image and fills in the details.
+The abstract from the paper is:
-The abstract of the paper is the following:
 *We present a neural network structure, ControlNet, to control pretrained large diffusion models to support additional input conditions. The ControlNet learns task-specific conditions in an end-to-end way, and the learning is robust even when the training dataset is small (< 50k). Moreover, training a ControlNet is as fast as fine-tuning a diffusion model, and the model can be trained on a personal devices. Alternatively, if powerful computation clusters are available, the model can scale to large amounts (millions to billions) of data. We report that large diffusion models like Stable Diffusion can be augmented with ControlNets to enable conditional inputs like edge maps, segmentation maps, keypoints, etc. This may enrich the methods to control large diffusion models and further facilitate related applications.*
-This model was contributed by the community contributor [takuma104](https://huggingface.co/takuma104) ❤️ .
+This model was contributed by [takuma104](https://huggingface.co/takuma104). ❤️
-Resources:
-* [Paper](https://arxiv.org/abs/2302.05543)
-* [Original Code](https://github.com/lllyasviel/ControlNet)
-## Available Pipelines:
-| Pipeline | Tasks | Demo
+The original codebase can be found at [lllyasviel/ControlNet](https://github.com/lllyasviel/ControlNet).
-|---|---|:---:|
-| [StableDiffusionControlNetPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/controlnet/pipeline_controlnet.py) | *Text-to-Image Generation with ControlNet Conditioning* | [Colab Example](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/controlnet.ipynb)
-| [StableDiffusionControlNetImg2ImgPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/controlnet/pipeline_controlnet_img2img.py) | *Image-to-Image Generation with ControlNet Conditioning* | 
-| [StableDiffusionControlNetInpaintPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_controlnet_inpaint.py) | *Inpainting Generation with ControlNet Conditioning* | 
 ## Usage example

--- a/docs/source/en/api/pipelines/cycle_diffusion.mdx
+++ b/docs/source/en/api/pipelines/cycle_diffusion.mdx
@@ -12,89 +12,22 @@ specific language governing permissions and limitations under the License.
 # Cycle Diffusion
-## Overview
+Cycle Diffusion is a text guided image-to-image generation model proposed in [Unifying Diffusion Models' Latent Space, with Applications to CycleDiffusion and Guidance](https://huggingface.co/papers/2210.05559) by Chen Henry Wu, Fernando De la Torre.
-Cycle Diffusion is a Text-Guided Image-to-Image Generation model proposed in [Unifying Diffusion Models' Latent Space, with Applications to CycleDiffusion and Guidance](https://arxiv.org/abs/2210.05559) by Chen Henry Wu, Fernando De la Torre.
+The abstract from the paper is:
-The abstract of the paper is the following:
 *Diffusion models have achieved unprecedented performance in generative modeling. The commonly-adopted formulation of the latent code of diffusion models is a sequence of gradually denoised samples, as opposed to the simpler (e.g., Gaussian) latent space of GANs, VAEs, and normalizing flows. This paper provides an alternative, Gaussian formulation of the latent space of various diffusion models, as well as an invertible DPM-Encoder that maps images into the latent space. While our formulation is purely based on the definition of diffusion models, we demonstrate several intriguing consequences. (1) Empirically, we observe that a common latent space emerges from two diffusion models trained independently on related domains. In light of this finding, we propose CycleDiffusion, which uses DPM-Encoder for unpaired image-to-image translation. Furthermore, applying CycleDiffusion to text-to-image diffusion models, we show that large-scale text-to-image diffusion models can be used as zero-shot image-to-image editors. (2) One can guide pre-trained diffusion models and GANs by controlling the latent codes in a unified, plug-and-play formulation based on energy-based models. Using the CLIP model and a face recognition model as guidance, we demonstrate that diffusion models have better coverage of low-density sub-populations and individuals than GANs.*
-*Tips*:
+<Tip>
- The Cycle Diffusion pipeline is fully compatible with any [Stable Diffusion](./stable_diffusion) checkpoints
- Currently Cycle Diffusion only works with the [`DDIMScheduler`].
-*Example*:
-In the following we should how to best use the [`CycleDiffusionPipeline`]
-```python
-import requests
-import torch
-from PIL import Image
-from io import BytesIO
-from diffusers import CycleDiffusionPipeline, DDIMScheduler
-# load the pipeline
-# make sure you're logged in with `huggingface-cli login`
-model_id_or_path = "CompVis/stable-diffusion-v1-4"
-scheduler = DDIMScheduler.from_pretrained(model_id_or_path, subfolder="scheduler")
-pipe = CycleDiffusionPipeline.from_pretrained(model_id_or_path, scheduler=scheduler).to("cuda")
-# let's download an initial image
-url = "https://raw.githubusercontent.com/ChenWu98/cycle-diffusion/main/data/dalle2/An%20astronaut%20riding%20a%20horse.png"
-response = requests.get(url)
-init_image = Image.open(BytesIO(response.content)).convert("RGB")
-init_image = init_image.resize((512, 512))
-init_image.save("horse.png")
-# let's specify a prompt
+Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
-source_prompt = "An astronaut riding a horse"
-prompt = "An astronaut riding an elephant"
-# call the pipeline
+</Tip>
-image = pipe(
-    prompt=prompt,
-    source_prompt=source_prompt,
-    image=init_image,
-    num_inference_steps=100,
-    eta=0.1,
-    strength=0.8,
-    guidance_scale=2,
-    source_guidance_scale=1,
-).images[0]
-image.save("horse_to_elephant.png")
-# let's try another example
-# See more samples at the original repo: https://github.com/ChenWu98/cycle-diffusion
-url = "https://raw.githubusercontent.com/ChenWu98/cycle-diffusion/main/data/dalle2/A%20black%20colored%20car.png"
-response = requests.get(url)
-init_image = Image.open(BytesIO(response.content)).convert("RGB")
-init_image = init_image.resize((512, 512))
-init_image.save("black.png")
-source_prompt = "A black colored car"
-prompt = "A blue colored car"
-# call the pipeline
-torch.manual_seed(0)
-image = pipe(
-    prompt=prompt,
-    source_prompt=source_prompt,
-    image=init_image,
-    num_inference_steps=100,
-    eta=0.1,
-    strength=0.85,
-    guidance_scale=3,
-    source_guidance_scale=1,
-).images[0]
-image.save("black_to_blue.png")
-```
 ## CycleDiffusionPipeline
 [[autodoc]] CycleDiffusionPipeline
 	- all
 	- __call__
+## StableDiffusionPiplineOutput
+[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput
\ No newline at end of file
--- a/docs/source/en/api/pipelines/dance_diffusion.mdx
+++ b/docs/source/en/api/pipelines/dance_diffusion.mdx
@@ -12,23 +12,22 @@ specific language governing permissions and limitations under the License.
 # Dance Diffusion
-## Overview
+[Dance Diffusion](https://github.com/Harmonai-org/sample-generator) is by Zach Evans.
-[Dance Diffusion](https://github.com/Harmonai-org/sample-generator) by Zach Evans.
+Dance Diffusion is the first in a suite of generative audio tools for producers and musicians released by [Harmonai](https://github.com/Harmonai-org).
-Dance Diffusion is the first in a suite of generative audio tools for producers and musicians to be released by Harmonai.
+The original codebase of this implementation can be found at [Harmonai-org](https://github.com/Harmonai-org/sample-generator).
-For more info or to get involved in the development of these tools, please visit https://harmonai.org and fill out the form on the front page.
-The original codebase of this implementation can be found [here](https://github.com/Harmonai-org/sample-generator).
+<Tip>
-## Available Pipelines:
+Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
-| Pipeline | Tasks | Colab
-|---|---|:---:|
-| [pipeline_dance_diffusion.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/dance_diffusion/pipeline_dance_diffusion.py) | *Unconditional Audio Generation* | - |
+</Tip>
 ## DanceDiffusionPipeline
 [[autodoc]] DanceDiffusionPipeline
 	- all
 	- __call__
+## AudioPipelineOutput
+[[autodoc]] pipelines.AudioPipelineOutput
\ No newline at end of file
--- a/docs/source/en/api/pipelines/ddim.mdx
+++ b/docs/source/en/api/pipelines/ddim.mdx
@@ -12,25 +12,18 @@ specific language governing permissions and limitations under the License.
 # DDIM
-## Overview
+[Denoising Diffusion Implicit Models](https://huggingface.co/papers/2010.02502) (DDIM) by Jiaming Song, Chenlin Meng and Stefano Ermon.
-[Denoising Diffusion Implicit Models](https://arxiv.org/abs/2010.02502) (DDIM) by Jiaming Song, Chenlin Meng and Stefano Ermon.
+The abstract from the paper is:
-The abstract of the paper is the following:
+*Denoising diffusion probabilistic models (DDPMs) have achieved high quality image generation without adversarial training, yet they require simulating a Markov chain for many steps to produce a sample. To accelerate sampling, we present denoising diffusion implicit models (DDIMs), a more efficient class of iterative implicit probabilistic models with the same training procedure as DDPMs. In DDPMs, the generative process is defined as the reverse of a Markovian diffusion process. We construct a class of non-Markovian diffusion processes that lead to the same training objective, but whose reverse process can be much faster to sample from. We empirically demonstrate that DDIMs can produce high quality samples 10× to 50× faster in terms of wall-clock time compared to DDPMs, allow us to trade off computation for sample quality, and can perform semantically meaningful image interpolation directly in the latent space.*
-Denoising diffusion probabilistic models (DDPMs) have achieved high quality image generation without adversarial training, yet they require simulating a Markov chain for many steps to produce a sample. To accelerate sampling, we present denoising diffusion implicit models (DDIMs), a more efficient class of iterative implicit probabilistic models with the same training procedure as DDPMs. In DDPMs, the generative process is defined as the reverse of a Markovian diffusion process. We construct a class of non-Markovian diffusion processes that lead to the same training objective, but whose reverse process can be much faster to sample from. We empirically demonstrate that DDIMs can produce high quality samples 10× to 50× faster in terms of wall-clock time compared to DDPMs, allow us to trade off computation for sample quality, and can perform semantically meaningful image interpolation directly in the latent space.
-The original codebase of this paper can be found here: [ermongroup/ddim](https://github.com/ermongroup/ddim).
-For questions, feel free to contact the author on [tsong.me](https://tsong.me/).
-## Available Pipelines:
-| Pipeline | Tasks | Colab
-|---|---|:---:|
-| [pipeline_ddim.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/ddim/pipeline_ddim.py) | *Unconditional Image Generation* | - |
+The original codebase can be found at [ermongroup/ddim](https://github.com/ermongroup/ddim).
 ## DDIMPipeline
 [[autodoc]] DDIMPipeline
 	- all
 	- __call__
+## ImagePipelineOutput
+[[autodoc]] pipelines.ImagePipelineOutput
\ No newline at end of file
--- a/docs/source/en/api/pipelines/ddpm.mdx
+++ b/docs/source/en/api/pipelines/ddpm.mdx
@@ -12,26 +12,24 @@ specific language governing permissions and limitations under the License.
 # DDPM
-## Overview
+[Denoising Diffusion Probabilistic Models](https://huggingface.co/papers/2006.11239) (DDPM) by Jonathan Ho, Ajay Jain and Pieter Abbeel proposes a diffusion based model of the same name. In the 🤗 Diffusers library, DDPM refers to the *discrete denoising scheduler* from the paper as well as the pipeline.
-[Denoising Diffusion Probabilistic Models](https://arxiv.org/abs/2006.11239) 
+The abstract from the paper is:
- (DDPM) by Jonathan Ho, Ajay Jain and Pieter Abbeel proposes the diffusion based model of the same name, but in the context of the 🤗 Diffusers library, DDPM refers to the discrete denoising scheduler from the paper as well as the pipeline.
-The abstract of the paper is the following:
+*We present high quality image synthesis results using diffusion probabilistic models, a class of latent variable models inspired by considerations from nonequilibrium thermodynamics. Our best results are obtained by training on a weighted variational bound designed according to a novel connection between diffusion probabilistic models and denoising score matching with Langevin dynamics, and our models naturally admit a progressive lossy decompression scheme that can be interpreted as a generalization of autoregressive decoding. On the unconditional CIFAR10 dataset, we obtain an Inception score of 9.46 and a state-of-the-art FID score of 3.17. On 256x256 LSUN, we obtain sample quality similar to ProgressiveGAN.*
-We present high quality image synthesis results using diffusion probabilistic models, a class of latent variable models inspired by considerations from nonequilibrium thermodynamics. Our best results are obtained by training on a weighted variational bound designed according to a novel connection between diffusion probabilistic models and denoising score matching with Langevin dynamics, and our models naturally admit a progressive lossy decompression scheme that can be interpreted as a generalization of autoregressive decoding. On the unconditional CIFAR10 dataset, we obtain an Inception score of 9.46 and a state-of-the-art FID score of 3.17. On 256x256 LSUN, we obtain sample quality similar to ProgressiveGAN.
+The original codebase can be found at [hohonathanho/diffusion](https://github.com/hojonathanho/diffusion).
-The original codebase of this paper can be found [here](https://github.com/hojonathanho/diffusion).
+<Tip>
+Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
-## Available Pipelines:
+</Tip>
-| Pipeline | Tasks | Colab
-|---|---|:---:|
-| [pipeline_ddpm.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/ddpm/pipeline_ddpm.py) | *Unconditional Image Generation* | - |
 # DDPMPipeline
 [[autodoc]] DDPMPipeline
 	- all
 	- __call__
+## ImagePipelineOutput
+[[autodoc]] pipelines.ImagePipelineOutput
--- a/docs/source/en/api/pipelines/diffedit.mdx
+++ b/docs/source/en/api/pipelines/diffedit.mdx
@@ -10,21 +10,17 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
 specific language governing permissions and limitations under the License.
 -->
-# Zero-shot Diffusion-based Semantic Image Editing with Mask Guidance
+# DiffEdit
-## Overview
+[DiffEdit: Diffusion-based semantic image editing with mask guidance](https://huggingface.co/papers/2210.11427) is by Guillaume Couairon, Jakob Verbeek, Holger Schwenk, and Matthieu Cord.
-[DiffEdit: Diffusion-based semantic image editing with mask guidance](https://arxiv.org/abs/2210.11427) by Guillaume Couairon, Jakob Verbeek, Holger Schwenk, and Matthieu Cord.
+The abstract from the paper is:
-The abstract of the paper is the following:
 *Image generation has recently seen tremendous advances, with diffusion models allowing to synthesize convincing images for a large variety of text prompts. In this article, we propose DiffEdit, a method to take advantage of text-conditioned diffusion models for the task of semantic image editing, where the goal is to edit an image based on a text query. Semantic image editing is an extension of image generation, with the additional constraint that the generated image should be as similar as possible to a given input image. Current editing methods based on diffusion models usually require to provide a mask, making the task much easier by treating it as a conditional inpainting task. In contrast, our main contribution is able to automatically generate a mask highlighting regions of the input image that need to be edited, by contrasting predictions of a diffusion model conditioned on different text prompts. Moreover, we rely on latent inference to preserve content in those regions of interest and show excellent synergies with mask-based diffusion. DiffEdit achieves state-of-the-art editing performance on ImageNet. In addition, we evaluate semantic image editing in more challenging settings, using images from the COCO dataset as well as text-based generated images.*
-Resources:
+The original codebase can be found at [Xiang-cd/DiffEdit-stable-diffusion](https://github.com/Xiang-cd/DiffEdit-stable-diffusion), and you can try it out in this [demo](https://blog.problemsolversguild.com/technical/research/2022/11/02/DiffEdit-Implementation.html).
-* [Paper](https://arxiv.org/abs/2210.11427).
+This pipeline was contributed by [clarencechen](https://github.com/clarencechen). ❤️
-* [Blog Post with Demo](https://blog.problemsolversguild.com/technical/research/2022/11/02/DiffEdit-Implementation.html).
-* [Implementation on Github](https://github.com/Xiang-cd/DiffEdit-stable-diffusion/).
 ## Tips 
@@ -51,14 +47,6 @@ below for more details.
    * Swap the `prompt` and `negative_prompt` in the arguments to call the pipeline to generate the final edited image.
 * Note that the source and target prompts, or their corresponding embeddings, can also be automatically generated. Please, refer to [this discussion](#generating-source-and-target-embeddings) for more details.
-## Available Pipelines:
-| Pipeline | Tasks
-|---|---|
-| [StableDiffusionDiffEditPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_diffedit.py) | *Text-Based Image Editing*
-<!-- TODO: add Colab -->
 ## Usage example
 ### Based on an input image with a caption
@@ -357,4 +345,4 @@ images[0].save("edited_image.png")
    - all
    - generate_mask
    - invert
-	- __call__
+    - __call__
\ No newline at end of file
--- a/docs/source/en/api/pipelines/dit.mdx
+++ b/docs/source/en/api/pipelines/dit.mdx
@@ -10,50 +10,26 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
 specific language governing permissions and limitations under the License.
 -->
-# Scalable Diffusion Models with Transformers (DiT)
+# DiT
-## Overview
+[Scalable Diffusion Models with Transformers](https://huggingface.co/papers/2212.09748) (DiT) is by William Peebles and Saining Xie.
-[Scalable Diffusion Models with Transformers](https://arxiv.org/abs/2212.09748) (DiT) by William Peebles and Saining Xie.
+The abstract from the paper is:
-The abstract of the paper is the following:
 *We explore a new class of diffusion models based on the transformer architecture. We train latent diffusion models of images, replacing the commonly-used U-Net backbone with a transformer that operates on latent patches. We analyze the scalability of our Diffusion Transformers (DiTs) through the lens of forward pass complexity as measured by Gflops. We find that DiTs with higher Gflops -- through increased transformer depth/width or increased number of input tokens -- consistently have lower FID. In addition to possessing good scalability properties, our largest DiT-XL/2 models outperform all prior diffusion models on the class-conditional ImageNet 512x512 and 256x256 benchmarks, achieving a state-of-the-art FID of 2.27 on the latter.*
-The original codebase of this paper can be found here: [facebookresearch/dit](https://github.com/facebookresearch/dit).
+The original codebase can be found at [facebookresearch/dit](https://github.com/facebookresearch/dit).
-## Available Pipelines:
-| Pipeline | Tasks | Colab
-|---|---|:---:|
-| [pipeline_dit.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/dit/pipeline_dit.py) | *Conditional Image Generation* | - |
-## Usage example
-```python
+<Tip>
-from diffusers import DiTPipeline, DPMSolverMultistepScheduler
-import torch
-pipe = DiTPipeline.from_pretrained("facebook/DiT-XL-2-256", torch_dtype=torch.float16)
+Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
-pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
-pipe = pipe.to("cuda")
-# pick words from Imagenet class labels
+</Tip>
-pipe.labels  # to print all available words
-# pick words that exist in ImageNet
-words = ["white shark", "umbrella"]
-class_ids = pipe.get_label_ids(words)
-generator = torch.manual_seed(33)
-output = pipe(class_labels=class_ids, num_inference_steps=25, generator=generator)
-image = output.images[0]  # label 'white shark'
-```
 ## DiTPipeline
 [[autodoc]] DiTPipeline
 	- all
 	- __call__
+## ImagePipelineOutput
+[[autodoc]] pipelines.ImagePipelineOutput
\ No newline at end of file
--- a/docs/source/en/api/pipelines/latent_diffusion.mdx
+++ b/docs/source/en/api/pipelines/latent_diffusion.mdx
@@ -12,31 +12,19 @@ specific language governing permissions and limitations under the License.
 # Latent Diffusion
-## Overview
+Latent Diffusion was proposed in [High-Resolution Image Synthesis with Latent Diffusion Models](https://huggingface.co/papers/2112.10752) by Robin Rombach, Andreas Blattmann, Dominik Lorenz, Patrick Esser, Björn Ommer.
-Latent Diffusion was proposed in [High-Resolution Image Synthesis with Latent Diffusion Models](https://arxiv.org/abs/2112.10752) by Robin Rombach, Andreas Blattmann, Dominik Lorenz, Patrick Esser, Björn Ommer.
+The abstract from the paper is:
-The abstract of the paper is the following:
 *By decomposing the image formation process into a sequential application of denoising autoencoders, diffusion models (DMs) achieve state-of-the-art synthesis results on image data and beyond. Additionally, their formulation allows for a guiding mechanism to control the image generation process without retraining. However, since these models typically operate directly in pixel space, optimization of powerful DMs often consumes hundreds of GPU days and inference is expensive due to sequential evaluations. To enable DM training on limited computational resources while retaining their quality and flexibility, we apply them in the latent space of powerful pretrained autoencoders. In contrast to previous work, training diffusion models on such a representation allows for the first time to reach a near-optimal point between complexity reduction and detail preservation, greatly boosting visual fidelity. By introducing cross-attention layers into the model architecture, we turn diffusion models into powerful and flexible generators for general conditioning inputs such as text or bounding boxes and high-resolution synthesis becomes possible in a convolutional manner. Our latent diffusion models (LDMs) achieve a new state of the art for image inpainting and highly competitive performance on various tasks, including unconditional image generation, semantic scene synthesis, and super-resolution, while significantly reducing computational requirements compared to pixel-based DMs.*
-The original codebase can be found [here](https://github.com/CompVis/latent-diffusion).
+The original codebase can be found at [Compvis/latent-diffusion](https://github.com/CompVis/latent-diffusion).
-## Tips:
- 
- 
- 
-## Available Pipelines:
+<Tip>
-| Pipeline | Tasks | Colab
+Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
-|---|---|:---:|
-| [pipeline_latent_diffusion.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/latent_diffusion/pipeline_latent_diffusion.py) | *Text-to-Image Generation* | - |
-| [pipeline_latent_diffusion_superresolution.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/latent_diffusion/pipeline_latent_diffusion_superresolution.py) | *Super Resolution* | - |
-## Examples:
+</Tip>
 ## LDMTextToImagePipeline
 [[autodoc]] LDMTextToImagePipeline
@@ -47,3 +35,6 @@ The original codebase can be found [here](https://github.com/CompVis/latent-diff
 [[autodoc]] LDMSuperResolutionPipeline
 	- all
 	- __call__
+## ImagePipelineOutput
+[[autodoc]] pipelines.ImagePipelineOutput
\ No newline at end of file
--- a/docs/source/en/api/pipelines/latent_diffusion_uncond.mdx
+++ b/docs/source/en/api/pipelines/latent_diffusion_uncond.mdx
@@ -12,31 +12,24 @@ specific language governing permissions and limitations under the License.
 # Unconditional Latent Diffusion
-## Overview
+Unconditional Latent Diffusion was proposed in [High-Resolution Image Synthesis with Latent Diffusion Models](https://huggingface.co/papers/2112.10752) by Robin Rombach, Andreas Blattmann, Dominik Lorenz, Patrick Esser, Björn Ommer.
-Unconditional Latent Diffusion was proposed in [High-Resolution Image Synthesis with Latent Diffusion Models](https://arxiv.org/abs/2112.10752) by Robin Rombach, Andreas Blattmann, Dominik Lorenz, Patrick Esser, Björn Ommer.
+The abstract from the paper is:
-The abstract of the paper is the following:
 *By decomposing the image formation process into a sequential application of denoising autoencoders, diffusion models (DMs) achieve state-of-the-art synthesis results on image data and beyond. Additionally, their formulation allows for a guiding mechanism to control the image generation process without retraining. However, since these models typically operate directly in pixel space, optimization of powerful DMs often consumes hundreds of GPU days and inference is expensive due to sequential evaluations. To enable DM training on limited computational resources while retaining their quality and flexibility, we apply them in the latent space of powerful pretrained autoencoders. In contrast to previous work, training diffusion models on such a representation allows for the first time to reach a near-optimal point between complexity reduction and detail preservation, greatly boosting visual fidelity. By introducing cross-attention layers into the model architecture, we turn diffusion models into powerful and flexible generators for general conditioning inputs such as text or bounding boxes and high-resolution synthesis becomes possible in a convolutional manner. Our latent diffusion models (LDMs) achieve a new state of the art for image inpainting and highly competitive performance on various tasks, including unconditional image generation, semantic scene synthesis, and super-resolution, while significantly reducing computational requirements compared to pixel-based DMs.*
-The original codebase can be found [here](https://github.com/CompVis/latent-diffusion).
+The original codebase can be found at [CompVis/latent-diffusion](https://github.com/CompVis/latent-diffusion).
-## Tips:
- 
+<Tip>
- 
- 
-## Available Pipelines:
+Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
-| Pipeline | Tasks | Colab
+</Tip>
-|---|---|:---:|
-| [pipeline_latent_diffusion_uncond.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/latent_diffusion_uncond/pipeline_latent_diffusion_uncond.py) | *Unconditional Image Generation* | - |
-## Examples:
 ## LDMPipeline
 [[autodoc]] LDMPipeline
 	- all
 	- __call__
+## ImagePipelineOutput
+[[autodoc]] pipelines.ImagePipelineOutput
--- a/docs/source/en/api/pipelines/model_editing.mdx
+++ b/docs/source/en/api/pipelines/model_editing.mdx
@@ -10,52 +10,26 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
 specific language governing permissions and limitations under the License.
 -->
-# Editing Implicit Assumptions in Text-to-Image Diffusion Models
+# Text-to-image model editing
-## Overview
+[Editing Implicit Assumptions in Text-to-Image Diffusion Models](https://huggingface.co/papers/2303.08084) is by Hadas Orgad, Bahjat Kawar, and Yonatan Belinkov. This pipeline enables editing diffusion model weights, such that its assumptions of a given concept are changed. The resulting change is expected to take effect in all prompt generations related to the edited concept.
-[Editing Implicit Assumptions in Text-to-Image Diffusion Models](https://arxiv.org/abs/2303.08084) by Hadas Orgad, Bahjat Kawar, and Yonatan Belinkov.
+The abstract from the paper is:
-The abstract of the paper is the following:
 *Text-to-image diffusion models often make implicit assumptions about the world when generating images. While some assumptions are useful (e.g., the sky is blue), they can also be outdated, incorrect, or reflective of social biases present in the training data. Thus, there is a need to control these assumptions without requiring explicit user input or costly re-training. In this work, we aim to edit a given implicit assumption in a pre-trained diffusion model. Our Text-to-Image Model Editing method, TIME for short, receives a pair of inputs: a "source" under-specified prompt for which the model makes an implicit assumption (e.g., "a pack of roses"), and a "destination" prompt that describes the same setting, but with a specified desired attribute (e.g., "a pack of blue roses"). TIME then updates the model's cross-attention layers, as these layers assign visual meaning to textual tokens. We edit the projection matrices in these layers such that the source prompt is projected close to the destination prompt. Our method is highly efficient, as it modifies a mere 2.2% of the model's parameters in under one second. To evaluate model editing approaches, we introduce TIMED (TIME Dataset), containing 147 source and destination prompt pairs from various domains. Our experiments (using Stable Diffusion) show that TIME is successful in model editing, generalizes well for related prompts unseen during editing, and imposes minimal effect on unrelated generations.*
-Resources:
+You can find additional information about model editing on the [project page](https://time-diffusion.github.io/), [original codebase](https://github.com/bahjat-kawar/time-diffusion), and try it out in a [demo](https://huggingface.co/spaces/bahjat-kawar/time-diffusion).
-* [Project Page](https://time-diffusion.github.io/).
-* [Paper](https://arxiv.org/abs/2303.08084).
-* [Original Code](https://github.com/bahjat-kawar/time-diffusion).
-* [Demo](https://huggingface.co/spaces/bahjat-kawar/time-diffusion).
-## Available Pipelines:
-| Pipeline | Tasks | Demo
-|---|---|:---:|
-| [StableDiffusionModelEditingPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_model_editing.py) | *Text-to-Image Model Editing* | [🤗 Space](https://huggingface.co/spaces/bahjat-kawar/time-diffusion)) |
-This pipeline enables editing the diffusion model weights, such that its assumptions on a given concept are changed. The resulting change is expected to take effect in all prompt generations pertaining to the edited concept.
-## Usage example
+<Tip>
-```python
+Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
-import torch
-from diffusers import StableDiffusionModelEditingPipeline
-model_ckpt = "CompVis/stable-diffusion-v1-4"
+</Tip>
-pipe = StableDiffusionModelEditingPipeline.from_pretrained(model_ckpt)
-pipe = pipe.to("cuda")
-source_prompt = "A pack of roses"
-destination_prompt = "A pack of blue roses"
-pipe.edit_model(source_prompt, destination_prompt)
-prompt = "A field of roses"
-image = pipe(prompt).images[0]
-image.save("field_of_roses.png")
-```
 ## StableDiffusionModelEditingPipeline
 [[autodoc]] StableDiffusionModelEditingPipeline
 	- __call__
 	- all
+## StableDiffusionPipelineOutput
+[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput
\ No newline at end of file
--- a/docs/source/en/api/pipelines/overview.mdx
+++ b/docs/source/en/api/pipelines/overview.mdx
@@ -12,108 +12,25 @@ specific language governing permissions and limitations under the License.
 # Pipelines
-Pipelines provide a simple way to run state-of-the-art diffusion models in inference.
+Pipelines provide a simple way to run state-of-the-art diffusion models in inference by bundling all of the necessary components (multiple independently-trained models, schedulers, and processors) into a single end-to-end class. Pipelines are flexible and they can be adapted to use different scheduler or even model components.
-Most diffusion systems consist of multiple independently-trained models and highly adaptable scheduler 
-components - all of which are needed to have a functioning end-to-end diffusion system.
-As an example, [Stable Diffusion](https://huggingface.co/blog/stable_diffusion) has three independently trained models:
+All pipelines are built from the base [`DiffusionPipeline`] class which provides basic functionality for loading, downloading, and saving all the components.
- [Autoencoder](./api/models#vae)
- [Conditional Unet](./api/models#UNet2DConditionModel)
- [CLIP text encoder](https://huggingface.co/docs/transformers/v4.27.1/en/model_doc/clip#transformers.CLIPTextModel)
- a scheduler component, [scheduler](./api/scheduler#pndm), 
- a [CLIPImageProcessor](https://huggingface.co/docs/transformers/v4.27.1/en/model_doc/clip#transformers.CLIPImageProcessor),
- as well as a [safety checker](./stable_diffusion#safety_checker).
-All of these components are necessary to run stable diffusion in inference even though they were trained 
-or created independently from each other.
-To that end, we strive to offer all open-sourced, state-of-the-art diffusion system under a unified API. 
+<Tip warning={true}>
-More specifically, we strive to provide pipelines that
- 1. can load the officially published weights and yield 1-to-1 the same outputs as the original implementation according to the corresponding paper (*e.g.* [LDMTextToImagePipeline](https://github.com/huggingface/diffusers/tree/main/src/diffusers/pipelines/latent_diffusion), uses the officially released weights of [High-Resolution Image Synthesis with Latent Diffusion Models](https://arxiv.org/abs/2112.10752)),
- 2. have a simple user interface to run the model in inference (see the [Pipelines API](#pipelines-api) section), 
- 3. are easy to understand with code that is self-explanatory and can be read along-side the official paper (see [Pipelines summary](#pipelines-summary)),
- 4. can easily be contributed by the community (see the [Contribution](#contribution) section).
-**Note** that pipelines do not (and should not) offer any training functionality. 
+Pipelines do not offer any training functionality. You'll notice PyTorch's autograd is disabled by decorating the [`~DiffusionPipeline.__call__`] method with a [`torch.no_grad`](https://pytorch.org/docs/stable/generated/torch.no_grad.html) decorator because pipelines should not be used for training. If you're interested in training, please take a look at the [Training](../traininig/overview) guides instead!
-If you are looking for *official* training examples, please have a look at [examples](https://github.com/huggingface/diffusers/tree/main/examples).
-## 🧨 Diffusers Summary
+</Tip>
-The following table summarizes all officially supported pipelines, their corresponding paper, and if 
+## DiffusionPipeline
-available a colab notebook to directly try them out.
+[[autodoc]] DiffusionPipeline
+	- all
+	- __call__
+	- device
+	- to
+	- components
-| Pipeline | Paper | Tasks | Colab
+## FlaxDiffusionPipeline
-|---|---|:---:|:---:|
-| [alt_diffusion](./alt_diffusion) | [**AltDiffusion**](https://arxiv.org/abs/2211.06679) | Image-to-Image Text-Guided Generation | -
-| [audio_diffusion](./audio_diffusion) | [**Audio Diffusion**](https://github.com/teticio/audio_diffusion.git) | Unconditional Audio Generation |
-| [controlnet](./api/pipelines/controlnet) | [**ControlNet with Stable Diffusion**](https://arxiv.org/abs/2302.05543) | Image-to-Image Text-Guided Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/controlnet.ipynb)
-| [cycle_diffusion](./cycle_diffusion) | [**Cycle Diffusion**](https://arxiv.org/abs/2210.05559) | Image-to-Image Text-Guided Generation |
-| [dance_diffusion](./dance_diffusion) | [**Dance Diffusion**](https://github.com/williamberman/diffusers.git) | Unconditional Audio Generation |
-| [ddpm](./ddpm) | [**Denoising Diffusion Probabilistic Models**](https://arxiv.org/abs/2006.11239) | Unconditional Image Generation |
-| [ddim](./ddim) | [**Denoising Diffusion Implicit Models**](https://arxiv.org/abs/2010.02502) | Unconditional Image Generation |
-| [if](./if) | [**IF**](https://github.com/deep-floyd/IF) | Image Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/deepfloyd_if_free_tier_google_colab.ipynb)
-| [if_img2img](./if) | [**IF**](https://github.com/deep-floyd/IF) | Image-to-Image Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/deepfloyd_if_free_tier_google_colab.ipynb)
-| [if_inpainting](./if) | [**IF**](https://github.com/deep-floyd/IF) | Image-to-Image Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/deepfloyd_if_free_tier_google_colab.ipynb)
-| [kandinsky](./kandinsky) | **Kandinsky** | Text-to-Image Generation | 
-| [kandinsky_inpaint](./kandinsky) | **Kandinsky** | Image-to-Image Generation | 
-| [kandinsky_img2img](./kandinsky) | **Kandinsksy** | Image-to-Image Generation |
-| [latent_diffusion](./latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Text-to-Image Generation | 
-| [latent_diffusion](./latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Super Resolution Image-to-Image | 
-| [latent_diffusion_uncond](./latent_diffusion_uncond) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752) | Unconditional Image Generation | 
-| [paint_by_example](./paint_by_example) | [**Paint by Example: Exemplar-based Image Editing with Diffusion Models**](https://arxiv.org/abs/2211.13227) | Image-Guided Image Inpainting | 
-| [paradigms](./paradigms) | [**Parallel Sampling of Diffusion Models**](https://arxiv.org/abs/2305.16317) | Text-to-Image Generation |
-| [pndm](./pndm) | [**Pseudo Numerical Methods for Diffusion Models on Manifolds**](https://arxiv.org/abs/2202.09778) | Unconditional Image Generation | 
-| [score_sde_ve](./score_sde_ve) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation | 
-| [score_sde_vp](./score_sde_vp) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation | 
-| [semantic_stable_diffusion](./semantic_stable_diffusion) | [**SEGA: Instructing Diffusion using Semantic Dimensions**](https://arxiv.org/abs/2301.12247) | Text-to-Image Generation |
-| [stable_diffusion_adapter](./stable_diffusion/adapter) | [**T2I-Adapter**](https://arxiv.org/abs/2302.08453) | Image-to-Image Text-Guided Generation with Adapters | -
-| [stable_diffusion_text2img](./stable_diffusion/text2img) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-to-Image Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/training_example.ipynb)
-| [stable_diffusion_img2img](./stable_diffusion/img2img) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Image-to-Image Text-Guided Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/image_2_image_using_diffusers.ipynb)
-| [stable_diffusion_inpaint](./stable_diffusion/inpaint) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-Guided Image Inpainting | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/in_painting_with_stable_diffusion_using_diffusers.ipynb)
-| [stable_diffusion_panorama](./stable_diffusion/panorama) | [**MultiDiffusion: Fusing Diffusion Paths for Controlled Image Generation**](https://arxiv.org/abs/2302.08113) | Text-Guided Panorama View Generation |
-| [stable_diffusion_pix2pix](./stable_diffusion/pix2pix) | [**InstructPix2Pix: Learning to Follow Image Editing Instructions**](https://arxiv.org/abs/2211.09800) | Text-Based Image Editing |
-| [stable_diffusion_pix2pix_zero](./stable_diffusion/pix2pix_zero) | [**Zero-shot Image-to-Image Translation**](https://arxiv.org/abs/2302.03027) | Text-Based Image Editing |
-| [stable_diffusion_attend_and_excite](./stable_diffusion/attend_and_excite) | [**Attend and Excite: Attention-Based Semantic Guidance for Text-to-Image Diffusion Models**](https://arxiv.org/abs/2301.13826) | Text-to-Image Generation |
-| [stable_diffusion_self_attention_guidance](./stable_diffusion/self_attention_guidance) | [**Self-Attention Guidance**](https://arxiv.org/abs/2210.00939) | Text-to-Image Generation |
-| [stable_diffusion_image_variation](./stable_diffusion/image_variation) | [**Stable Diffusion Image Variations**](https://github.com/LambdaLabsML/lambda-diffusers#stable-diffusion-image-variations) | Image-to-Image Generation |
-| [stable_diffusion_latent_upscale](./stable_diffusion/latent_upscale) | [**Stable Diffusion Latent Upscaler**](https://twitter.com/StabilityAI/status/1590531958815064065) | Text-Guided Super Resolution Image-to-Image |
-| [stable_diffusion_2](./stable_diffusion/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Image Inpainting | 
-| [stable_diffusion_2](./stable_diffusion/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Depth-to-Image Text-Guided Generation |
-| [stable_diffusion_2](./stable_diffusion/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Super Resolution Image-to-Image |
-| [stable_diffusion_safe](./stable_diffusion_safe) | [**Safe Stable Diffusion**](https://arxiv.org/abs/2211.05105) | Text-Guided Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/ml-research/safe-latent-diffusion/blob/main/examples/Safe%20Latent%20Diffusion.ipynb)
-| [stable_unclip](./stable_unclip) | **Stable unCLIP** | Text-to-Image Generation |
-| [stable_unclip](./stable_unclip) | **Stable unCLIP** | Image-to-Image Text-Guided Generation |
-| [stochastic_karras_ve](./stochastic_karras_ve) | [**Elucidating the Design Space of Diffusion-Based Generative Models**](https://arxiv.org/abs/2206.00364) | Unconditional Image Generation |
-| [text_to_video_sd](./api/pipelines/text_to_video) | [**Modelscope's Text-to-video-synthesis Model in Open Domain**](https://modelscope.cn/models/damo/text-to-video-synthesis/summary) | Text-to-Video Generation |
-| [unclip](./unclip) | [**Hierarchical Text-Conditional Image Generation with CLIP Latents](https://arxiv.org/abs/2204.06125) | Text-to-Image Generation |
-| [versatile_diffusion](./versatile_diffusion) | [**Versatile Diffusion: Text, Images and Variations All in One Diffusion Model**](https://arxiv.org/abs/2211.08332) | Text-to-Image Generation | 
-| [versatile_diffusion](./versatile_diffusion) | [**Versatile Diffusion: Text, Images and Variations All in One Diffusion Model**](https://arxiv.org/abs/2211.08332) | Image Variations Generation | 
-| [versatile_diffusion](./versatile_diffusion) | [**Versatile Diffusion: Text, Images and Variations All in One Diffusion Model**](https://arxiv.org/abs/2211.08332) | Dual Image and Text Guided Generation | 
-| [vq_diffusion](./vq_diffusion) | [**Vector Quantized Diffusion Model for Text-to-Image Synthesis**](https://arxiv.org/abs/2111.14822) | Text-to-Image Generation | 
-| [text_to_video_zero](./text_to_video_zero) | [**Text2Video-Zero: Text-to-Image Diffusion Models are Zero-Shot Video Generators**](https://arxiv.org/abs/2303.13439) | Text-to-Video Generation |
+[[autodoc]] pipelines.pipeline_flax_utils.FlaxDiffusionPipeline
-**Note**: Pipelines are simple examples of how to play around with the diffusion systems as described in the corresponding papers. 
-However, most of them can be adapted to use different scheduler components or even different model components. Some pipeline examples are shown in the [Examples](#examples) below.
-## Pipelines API
-Diffusion models often consist of multiple independently-trained models or other previously existing components. 
-Each model has been trained independently on a different task and the scheduler can easily be swapped out and replaced with a different one. 
-During inference, we however want to be able to easily load all components and use them in inference - even if one component, *e.g.* CLIP's text encoder, originates from a different library, such as [Transformers](https://github.com/huggingface/transformers). To that end, all pipelines provide the following functionality:
- [`from_pretrained` method](../diffusion_pipeline) that accepts a Hugging Face Hub repository id, *e.g.* [runwayml/stable-diffusion-v1-5](https://huggingface.co/runwayml/stable-diffusion-v1-5) or a path to a local directory, *e.g.*
-"./stable-diffusion". To correctly retrieve which models and components should be loaded, one has to provide a `model_index.json` file, *e.g.* [runwayml/stable-diffusion-v1-5/model_index.json](https://huggingface.co/runwayml/stable-diffusion-v1-5/blob/main/model_index.json), which defines all components that should be
-loaded into the pipelines. More specifically, for each model/component one needs to define the format `<name>: ["<library>", "<class name>"]`. `<name>` is the attribute name given to the loaded instance of `<class name>` which can be found in the library or pipeline folder called `"<library>"`.
- [`save_pretrained`](../diffusion_pipeline) that accepts a local path, *e.g.* `./stable-diffusion` under which all models/components of the pipeline will be saved. For each component/model a folder is created inside the local path that is named after the given attribute name, *e.g.* `./stable_diffusion/unet`. 
-In addition, a `model_index.json` file is created at the root of the local path, *e.g.* `./stable_diffusion/model_index.json` so that the complete pipeline can again be instantiated 
-from the local path.
- [`to`](../diffusion_pipeline) which accepts a `string` or `torch.device` to move all models that are of type `torch.nn.Module` to the passed device. The behavior is fully analogous to [PyTorch's `to` method](https://pytorch.org/docs/stable/generated/torch.nn.Module.html#torch.nn.Module.to).
- [`__call__`] method to use the pipeline in inference. `__call__` defines inference logic of the pipeline and should ideally encompass all aspects of it, from pre-processing to forwarding tensors to the different models and schedulers, as well as post-processing. The API of the `__call__` method can strongly vary from pipeline to pipeline. *E.g.* a text-to-image pipeline, such as [`StableDiffusionPipeline`](./stable_diffusion) should accept among other things the text prompt to generate the image. A pure image generation pipeline, such as [DDPMPipeline](https://github.com/huggingface/diffusers/tree/main/src/diffusers/pipelines/ddpm) on the other hand can be run without providing any inputs. To better understand what inputs can be adapted for 
-each pipeline, one should look directly into the respective pipeline.
-**Note**: All pipelines have PyTorch's autograd disabled by decorating the `__call__` method with a [`torch.no_grad`](https://pytorch.org/docs/stable/generated/torch.no_grad.html) decorator because pipelines should
-not be used for training. If you want to store the gradients during the forward pass, we recommend writing your own pipeline, see also our [community-examples](https://github.com/huggingface/diffusers/tree/main/examples/community).
--- a/docs/source/en/api/pipelines/paint_by_example.mdx
+++ b/docs/source/en/api/pipelines/paint_by_example.mdx
@@ -12,63 +12,28 @@ specific language governing permissions and limitations under the License.
 # PaintByExample
-## Overview
+[Paint by Example: Exemplar-based Image Editing with Diffusion Models](https://huggingface.co/papers/2211.13227) is by Binxin Yang, Shuyang Gu, Bo Zhang, Ting Zhang, Xuejin Chen, Xiaoyan Sun, Dong Chen, Fang Wen.
-[Paint by Example: Exemplar-based Image Editing with Diffusion Models](https://arxiv.org/abs/2211.13227) by Binxin Yang, Shuyang Gu, Bo Zhang, Ting Zhang, Xuejin Chen, Xiaoyan Sun, Dong Chen, Fang Wen.
+The abstract from the paper is:
-The abstract of the paper is the following:
 *Language-guided image editing has achieved great success recently. In this paper, for the first time, we investigate exemplar-guided image editing for more precise control. We achieve this goal by leveraging self-supervised training to disentangle and re-organize the source image and the exemplar. However, the naive approach will cause obvious fusing artifacts. We carefully analyze it and propose an information bottleneck and strong augmentations to avoid the trivial solution of directly copying and pasting the exemplar image. Meanwhile, to ensure the controllability of the editing process, we design an arbitrary shape mask for the exemplar image and leverage the classifier-free guidance to increase the similarity to the exemplar image. The whole framework involves a single forward of the diffusion model without any iterative optimization. We demonstrate that our method achieves an impressive performance and enables controllable editing on in-the-wild images with high fidelity.*
-The original codebase can be found [here](https://github.com/Fantasy-Studio/Paint-by-Example).
+The original codebase can be found at [Fantasy-Studio/Paint-by-Example](https://github.com/Fantasy-Studio/Paint-by-Example), and you can try it out in a [demo](https://huggingface.co/spaces/Fantasy-Studio/Paint-by-Example).
-## Available Pipelines:
-| Pipeline | Tasks | Colab
-|---|---|:---:|
-| [pipeline_paint_by_example.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/paint_by_example/pipeline_paint_by_example.py) | *Image-Guided Image Painting* | - |
 ## Tips
- PaintByExample is supported by the official [Fantasy-Studio/Paint-by-Example](https://huggingface.co/Fantasy-Studio/Paint-by-Example) checkpoint. The checkpoint has been warm-started from the [CompVis/stable-diffusion-v1-4](https://huggingface.co/CompVis/stable-diffusion-v1-4) and with the objective to inpaint partly masked images conditioned on example / reference images
+PaintByExample is supported by the official [Fantasy-Studio/Paint-by-Example](https://huggingface.co/Fantasy-Studio/Paint-by-Example) checkpoint. The checkpoint is warm-started from [CompVis/stable-diffusion-v1-4](https://huggingface.co/CompVis/stable-diffusion-v1-4) to inpaint partly masked images conditioned on example and reference images.
- To quickly demo *PaintByExample*, please have a look at [this demo](https://huggingface.co/spaces/Fantasy-Studio/Paint-by-Example)
- You can run the following code snippet as an example:
-```python
-# !pip install diffusers transformers
-import PIL
-import requests
-import torch
-from io import BytesIO
-from diffusers import DiffusionPipeline
+<Tip>
-def download_image(url):
+Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
-    response = requests.get(url)
-    return PIL.Image.open(BytesIO(response.content)).convert("RGB")
+</Tip>
-img_url = "https://raw.githubusercontent.com/Fantasy-Studio/Paint-by-Example/main/examples/image/example_1.png"
-mask_url = "https://raw.githubusercontent.com/Fantasy-Studio/Paint-by-Example/main/examples/mask/example_1.png"
-example_url = "https://raw.githubusercontent.com/Fantasy-Studio/Paint-by-Example/main/examples/reference/example_1.jpg"
-init_image = download_image(img_url).resize((512, 512))
-mask_image = download_image(mask_url).resize((512, 512))
-example_image = download_image(example_url).resize((512, 512))
-pipe = DiffusionPipeline.from_pretrained(
-    "Fantasy-Studio/Paint-by-Example",
-    torch_dtype=torch.float16,
-)
-pipe = pipe.to("cuda")
-image = pipe(image=init_image, mask_image=mask_image, example_image=example_image).images[0]
-image
-```
 ## PaintByExamplePipeline
 [[autodoc]] PaintByExamplePipeline
 	- all
 	- __call__
+## StableDiffusionPipelineOutput
+[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput
\ No newline at end of file
--- a/docs/source/en/api/pipelines/panorama.mdx
+++ b/docs/source/en/api/pipelines/panorama.mdx
@@ -10,76 +10,48 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
 specific language governing permissions and limitations under the License.
 -->
-# MultiDiffusion: Fusing Diffusion Paths for Controlled Image Generation
+# MultiDiffusion
-## Overview
+[MultiDiffusion: Fusing Diffusion Paths for Controlled Image Generation](https://huggingface.co/papers/2302.08113) is by Omer Bar-Tal, Lior Yariv, Yaron Lipman, and Tali Dekel.
-[MultiDiffusion: Fusing Diffusion Paths for Controlled Image Generation](https://arxiv.org/abs/2302.08113) by Omer Bar-Tal, Lior Yariv, Yaron Lipman, and Tali Dekel.
+The abstract from the paper is:
-The abstract of the paper is the following:
+*Recent advances in text-to-image generation with diffusion models present transformative capabilities in image quality. However, user controllability of the generated image, and fast adaptation to new tasks still remains an open challenge, currently mostly addressed by costly and long re-training and fine-tuning or ad-hoc adaptations to specific image generation tasks. In this work, we present MultiDiffusion, a unified framework that enables versatile and controllable image generation, using a pre-trained text-to-image diffusion model, without any further training or finetuning. At the center of our approach is a new generation process, based on an optimization task that binds together multiple diffusion generation processes with a shared set of parameters or constraints. We show that MultiDiffusion can be readily applied to generate high quality and diverse images that adhere to user-provided controls, such as desired aspect ratio (e.g., panorama), and spatial guiding signals, ranging from tight segmentation masks to bounding boxes.*
-*Recent advances in text-to-image generation with diffusion models present transformative capabilities in image quality. However, user controllability of the generated image, and fast adaptation to new tasks still remains an open challenge, currently mostly addressed by costly and long re-training and fine-tuning or ad-hoc adaptations to specific image generation tasks. In this work, we present MultiDiffusion, a unified framework that enables versatile and controllable image generation, using a pre-trained text-to-image diffusion model, without any further training or finetuning. At the center of our approach is a new generation process, based on an optimization task that binds together multiple diffusion generation processes with a shared set of parameters or constraints. We show that MultiDiffusion can be readily applied to generate high quality and diverse images that adhere to user-provided controls, such as desired aspect ratio (e.g., panorama), and spatial guiding signals, ranging from tight segmentation masks to bounding boxes.
+You can find additional information about MultiDiffusion on the [project page](https://multidiffusion.github.io/), [original codebase](https://github.com/omerbt/MultiDiffusion), and try it out in a [demo](https://huggingface.co/spaces/weizmannscience/MultiDiffusion).
-Resources:
+## Tips
-* [Project Page](https://multidiffusion.github.io/).
+While calling [`StableDiffusionPanoramaPipeline`], it's possible to specify the `view_batch_size` parameter to be > 1. 
-* [Paper](https://arxiv.org/abs/2302.08113).
+For some GPUs with high performance, this can speedup the generation process and increase VRAM usage.
-* [Original Code](https://github.com/omerbt/MultiDiffusion).
-* [Demo](https://huggingface.co/spaces/weizmannscience/MultiDiffusion).
-## Available Pipelines:
+To generate panorama-like images make sure you pass the width parameter accordingly. We recommend a width value of 2048 which is the default.
-| Pipeline | Tasks | Demo
-|---|---|:---:|
-| [StableDiffusionPanoramaPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_panorama.py) | *Text-Guided Panorama View Generation* | [🤗 Space](https://huggingface.co/spaces/weizmannscience/MultiDiffusion)) |
-<!-- TODO: add Colab -->
-## Usage example
-```python
-import torch
-from diffusers import StableDiffusionPanoramaPipeline, DDIMScheduler
-model_ckpt = "stabilityai/stable-diffusion-2-base"
-scheduler = DDIMScheduler.from_pretrained(model_ckpt, subfolder="scheduler")
-pipe = StableDiffusionPanoramaPipeline.from_pretrained(model_ckpt, scheduler=scheduler, torch_dtype=torch.float16)
-pipe = pipe.to("cuda")
-prompt = "a photo of the dolomites"
-image = pipe(prompt).images[0]
-image.save("dolomites.png")
-```
-<Tip>
-While calling this pipeline, it's possible to specify the `view_batch_size` to have a >1 value. 
-For some GPUs with high performance, higher a `view_batch_size`, can speedup the generation
-and increase the VRAM usage.
-</Tip>
-<Tip>
 Circular padding is applied to ensure there are no stitching artifacts when working with 
-panoramas that needs to seamlessly transition from the rightmost part to the leftmost part. 
+panoramas to ensure a seamless transition from the rightmost part to the leftmost part. 
 By enabling circular padding (set `circular_padding=True`), the operation applies additional 
 crops after the rightmost point of the image, allowing the model to "see” the transition 
 from the rightmost part to the leftmost part. This helps maintain visual consistency in 
 a 360-degree sense and creates a proper “panorama” that can be viewed using 360-degree 
-panorama viewers. When decoding latents in StableDiffusion, circular padding is applied 
+panorama viewers. When decoding latents in Stable Diffusion, circular padding is applied 
 to ensure that the decoded latents match in the RGB space.
-Without circular padding, there is a stitching artifact (default):
+For example, without circular padding, there is a stitching artifact (default):
 ![img](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/indoor_%20no_circular_padding.png)
-With circular padding, the right and the left parts are matching (`circular_padding=True`):
+But with circular padding, the right and the left parts are matching (`circular_padding=True`):
 ![img](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/indoor_%20circular_padding.png)
+<Tip>
+Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
 </Tip>
 ## StableDiffusionPanoramaPipeline
 [[autodoc]] StableDiffusionPanoramaPipeline
 	- __call__
 	- all
+## StableDiffusionPipelineOutput
+[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput
\ No newline at end of file
--- a/docs/source/en/api/pipelines/paradigms.mdx
+++ b/docs/source/en/api/pipelines/paradigms.mdx
@@ -10,74 +10,45 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
 specific language governing permissions and limitations under the License.
 -->
-# Parallel Sampling of Diffusion Models (ParaDiGMS)
+# Parallel Sampling of Diffusion Models
-## Overview
+[Parallel Sampling of Diffusion Models](https://huggingface.co/papers/2305.16317) is by Andy Shih, Suneel Belkhale, Stefano Ermon, Dorsa Sadigh, Nima Anari.
-[Parallel Sampling of Diffusion Models](https://arxiv.org/abs/2305.16317) by Andy Shih, Suneel Belkhale, Stefano Ermon, Dorsa Sadigh, Nima Anari.
+The abstract from the paper is:
-The abstract of the paper is the following:
 *Diffusion models are powerful generative models but suffer from slow sampling, often taking 1000 sequential denoising steps for one sample. As a result, considerable efforts have been directed toward reducing the number of denoising steps, but these methods hurt sample quality. Instead of reducing the number of denoising steps (trading quality for speed), in this paper we explore an orthogonal approach: can we run the denoising steps in parallel (trading compute for speed)? In spite of the sequential nature of the denoising steps, we show that surprisingly it is possible to parallelize sampling via Picard iterations, by guessing the solution of future denoising steps and iteratively refining until convergence. With this insight, we present ParaDiGMS, a novel method to accelerate the sampling of pretrained diffusion models by denoising multiple steps in parallel. ParaDiGMS is the first diffusion sampling method that enables trading compute for speed and is even compatible with existing fast sampling techniques such as DDIM and DPMSolver. Using ParaDiGMS, we improve sampling speed by 2-4x across a range of robotics and image generation models, giving state-of-the-art sampling speeds of 0.2s on 100-step DiffusionPolicy and 16s on 1000-step StableDiffusion-v2 with no measurable degradation of task reward, FID score, or CLIP score.*
-Resources:
+The original codebase can be found at [AndyShih12/paradigms](https://github.com/AndyShih12/paradigms), and the pipeline was contributed by [AndyShih12](https://github.com/AndyShih12). ❤️
-* [Paper](https://arxiv.org/abs/2305.16317).
-* [Original Code](https://github.com/AndyShih12/paradigms).
-## Available Pipelines:
-| Pipeline | Tasks | Demo
-|---|---|:---:|
-| [StableDiffusionParadigmsPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_paradigms.py) | *Faster Text-to-Image Generation* | |
-This pipeline was contributed by [`AndyShih12`](https://github.com/AndyShih12) in this [PR](https://github.com/huggingface/diffusers/pull/3716/).
-## Usage example
-```python
-import torch
-from diffusers import DDPMParallelScheduler
-from diffusers import StableDiffusionParadigmsPipeline
-scheduler = DDPMParallelScheduler.from_pretrained("runwayml/stable-diffusion-v1-5", subfolder="scheduler")
-pipe = StableDiffusionParadigmsPipeline.from_pretrained(
+## Tips
-    "runwayml/stable-diffusion-v1-5", scheduler=scheduler, torch_dtype=torch.float16
-)
-pipe = pipe.to("cuda")
-ngpu, batch_per_device = torch.cuda.device_count(), 5
-pipe.wrapped_unet = torch.nn.DataParallel(pipe.unet, device_ids=[d for d in range(ngpu)])
-prompt = "a photo of an astronaut riding a horse on mars"
-image = pipe(prompt, parallel=ngpu * batch_per_device, num_inference_steps=1000).images[0]
-```
-<Tip>
 This pipeline improves sampling speed by running denoising steps in parallel, at the cost of increased total FLOPs.
 Therefore, it is better to call this pipeline when running on multiple GPUs. Otherwise, without enough GPU bandwidth
 sampling may be even slower than sequential sampling.
 The two parameters to play with are `parallel` (batch size) and `tolerance`. 
- If it fits in memory, for 1000-step DDPM you can aim for a batch size of around 100 
+- If it fits in memory, for a 1000-step DDPM you can aim for a batch size of around 100 
-(e.g. 8 GPUs and batch_per_device=12 to get parallel=96). Higher batch size
+(for example, 8 GPUs and `batch_per_device=12` to get `parallel=96`). A higher batch size
 may not fit in memory, and lower batch size gives less parallelism. 
 - For tolerance, using a higher tolerance may get better speedups but can risk sample quality degradation. 
-If there is quality degradation with the default tolerance, then use a lower tolerance (e.g. 0.001).
+If there is quality degradation with the default tolerance, then use a lower tolerance like `0.001`.
-For 1000-step DDPM on 8 A100 GPUs, you can expect around a 3x speedup by StableDiffusionParadigmsPipeline instead of StableDiffusionPipeline
+For a 1000-step DDPM on 8 A100 GPUs, you can expect around a 3x speedup from [`StableDiffusionParadigmsPipeline`] compared to the [`StableDiffusionPipeline`]
-by setting parallel=80 and tolerance=0.1.
+by setting `parallel=80` and `tolerance=0.1`.
-</Tip>
+🤗 Diffusers offers [distributed inference support](../training/distributed_inference) for generating multiple prompts
+in parallel on multiple GPUs. But [`StableDiffusionParadigmsPipeline`] is designed for speeding up sampling of a single prompt by using multiple GPUs.
 <Tip>
-Diffusers also offers distributed inference support for generating multiple prompts
-in parallel on multiple GPUs. Check out the docs [here](https://huggingface.co/docs/diffusers/main/en/training/distributed_inference).
-In contrast, this pipeline is designed for speeding up sampling of a single prompt (by using multiple GPUs).
+Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
 </Tip>
 ## StableDiffusionParadigmsPipeline
 [[autodoc]] StableDiffusionParadigmsPipeline
 	- __call__
 	- all
+## StableDiffusionPipelineOutput
+[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput