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# Models
# Diffusion Pipeline
Diffusers contains pretrained models for popular algorithms and modules for creating the next set of diffusion models.
The primary function of these models is to denoise an input sample, by modeling the distribution $p_\theta(\mathbf{x}_{t-1}|\mathbf{x}_t)$.
The models are built on the base class ['ModelMixin'] that is a `torch.nn.module` with basic functionality for saving and loading models both locally and from the HuggingFace hub.
The [`DiffusionPipeline`] is the easiest way to load any pretrained diffusion pipeline from the [Hub](https://huggingface.co/models?library=diffusers) and to use it in inference.
## API
<Tip>
Models should provide the `def forward` function and initialization of the model.
All saving, loading, and utilities should be in the base ['ModelMixin'] class.
One should not use the Diffusion Pipeline class for training or fine-tuning a diffusion model. Individual
components of diffusion pipelines are usually trained individually, so we suggest to directly work
with [`UNetModel`] and [`UNetConditionModel`].
## Examples
</Tip>
- The ['UNetModel'] was proposed in [TODO](https://arxiv.org/) and has been used in paper1, paper2, paper3.
- Extensions of the ['UNetModel'] include the ['UNetGlideModel'] that uses attention and timestep embeddings for the [GLIDE](https://arxiv.org/abs/2112.10741) paper, the ['UNetGradTTS'] model from this [paper](https://arxiv.org/abs/2105.06337) for text-to-speech, ['UNetLDMModel'] for latent-diffusion models in this [paper](https://arxiv.org/abs/2112.10752), and the ['TemporalUNet'] used for time-series prediciton in this reinforcement learning [paper](https://arxiv.org/abs/2205.09991).
- TODO: mention VAE / SDE score estimation
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Any diffusion pipeline that is loaded with [`~DiffusionPipeline.from_pretrained`] will automatically
detect the pipeline type, *e.g.* [`StableDiffusionPipeline`] and consequently load each component of the
pipeline and pass them into the `__init__` function of the pipeline, *e.g.* [`~StableDiffusionPipeline.__init__`].
Any pipeline object can be saved locally with [`~DiffusionPipeline.save_pretrained`].
Stable Diffusion is a text-to-image _latent diffusion_ model created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/) and [LAION](https://laion.ai/). It's trained on 512x512 images from a subset of the [LAION-5B](https://laion.ai/blog/laion-5b/) dataset. This model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts. With its 860M UNet and 123M text encoder, the model is relatively lightweight and can run on consumer GPUs.
Latent diffusion is the research on top of which Stable Diffusion was built. It was proposed in [High-Resolution Image Synthesis with Latent Diffusion Models](https://arxiv.org/abs/2112.10752) by Robin Rombach, Andreas Blattmann, Dominik Lorenz, Patrick Esser, Björn Ommer. You can learn more details about it in the [specific pipeline for latent diffusion](pipelines/latent_diffusion) that is part of 🤗 Diffusers.
For more details about how Stable Diffusion works and how it differs from the base latent diffusion model, please refer to the official [launch announcement post](https://stability.ai/blog/stable-diffusion-announcement) and [this section of our own blog post](https://huggingface.co/blog/stable_diffusion#how-does-stable-diffusion-work).
## Tips
*Tips*:
- To tweak your prompts on a specific result you liked, you can generate your own latents, as demonstrated in the following notebook: [](https://colab.research.google.com/github/pcuenca/diffusers-examples/blob/main/notebooks/stable-diffusion-seeds.ipynb)